For the rest of this guide, we'll either use the generic Stable Diffusion v1. Unprompted is a highly modular extension for AUTOMATIC1111's Stable Diffusion Web UI that allows you to include various shortcodes in your prompts. To start using ChatGPT, go to chat. The model files used in the inference should be uploaded to the cloud before generate, which can be referred to the introduction of chapter Cloud Assets Management. Important: An Nvidia GPU with at least 10 GB is recommended. The company claims this is the fastest-ever local deployment of the tool on a smartphone. Hosted on Banana 🍌. Stable Diffusion. 6. 恭喜你发现了宝藏新博主🎉萌新的第一次投稿,望大家多多支持和关注保姆级stable diffusion + mov2mov 一键出ai视频做视频好累啊,视频做了一天,写扩展用了一天使用规约:请自行解决视频来源的授权问题，任何由于使用非授权视频进行转换造成的问题，需自行承担全部责任和一切后果，于mov2mov无关！任何. 指定した画像に近づくように画像生成する機能です。通常のプロンプトによる生成指定に加えて、追加でVGG16の特徴量を取得し、生成中の画像が指定したガイド画像に近づくよう、生成される画像をコントロールします。2. Hosted on Banana 🍌. Go to the bottom of the generation parameters and select the script. 手順3：PowerShellでコマンドを打ち込み、環境を構築する. Prompt: the description of the image the AI is going to generate. See the SDXL guide for an alternative setup with SD. 4/5 generated image and get the prompt to replicate that image/style. Contents. I used two different yet similar prompts and did 4 A/B studies with each prompt. Next, copy your API token and authenticate by setting it as an environment variable: export REPLICATE_API_TOKEN=<paste-your-token-here>. Set the batch size to 4 so that you can. It really depends on what you're using to run the Stable Diffusion. generating img2txt with the new v2. . 3 - One Step Closer to Reality Research Model - How to Build Protogen Running on Apple Silicon devices ? Try this instead. 5 model. It was pre-trained being conditioned on the ImageNet-1k classes. 08:41. Put this in the prompt text box. nsfw. It allows the model to generate contextualized images of the subject in different scenes, poses, and views. If you’ve saved new models in there while A1111 is running you can hit the blue refresh button to the right of the drop. The backbone. NAI is a model created by the company NovelAI modifying the Stable Diffusion architecture and training method. More awesome work from Christian Cantrell in his free plugin. Under the Generate button there is an Interrogate CLIP which when clicked will download the CLIP for reasoning about the Prompt of the image in the current image box and filling it to the prompt. A random selection of images created using AI text to image generator Stable Diffusion. Iterate if necessary: If the results are not satisfactory, adjust the filter parameters or try a different filter. . Aug 26, 2022. Whilst the then popular Waifu Diffusion was trained on SD + 300k anime images, NAI was trained on millions. Stable diffusionのイカしたテクニック、txt2imghdの仕組みを解説します。 簡単に試すことのできるGoogle Colabも添付しましたので、是非お試しください。 ↓の画像は、通常のtxt2imgとtxt2imghdで生成した画像を拡大して並べたものです。明らかに綺麗になっていること. Download: Installation: Extract anywhere (not a protected folder - NOT Program Files - preferrably a short custom path like D:/Apps/AI/), run StableDiffusionGui. safetensors (5. - use img2txt to generate the prompt and img2img to provide the starting point. novelai用了下，故意挑了些涩图tag，效果还可以 基于stable diffusion，操作和sd类似 他们的介绍文档 价格主要是订阅那一下有点贵，要10刀，送1000token 一张图5token(512*768)，细化什么的额外消耗token 这方面倒还好，就是买算力了… 充值token 10刀10000左右，其实还行Model Type. fix）を使っている方もいるかもしれません。 ですが、ハイレゾは大容量のVRAMが必要で、途中でエラーになって停止してしまうことがありま. Thanks to the passionate community, most new features come to this free Stable Diffusion GUI first. The StableDiffusionPipeline is capable of generating photorealistic images given any text input. pixray / text2image. Height. Install the Node. This guide will show you how to finetune the CompVis/stable-diffusion-v1-4 model on your own dataset with PyTorch and Flax. 5, ControlNet Linear/OpenPose, DeFlicker Resolve. Checkpoints (. Enjoy . Hi, yes you can mix two even more images with stable diffusion. ago. Contents. I originally tried this with DALL-E with similar prompts and the results are less appetizing. Stable Difussion Web UIのHires. The VD-basic is an image variation model with a single-flow. Run time and cost. exe"kaggle competitions download -c stable-diffusion-image-to-prompts unzip stable-diffusion-image-to-prompts. You can receive up to four options per prompt. By default, 🤗 Diffusers automatically loads these . for examples:"Logo of a pirate","logo of a sunglass with girl" or something complex like "logo of a ice-cream with snake" etc. Scroll to the bottom of the notebook to the Prompts section near the very bottom of the notebook. $0. Usually, higher is better but to a certain degree. AI不仅能够自动用文字生成画面，还能够对制定的图片扩展画面意外的内容，也就是根据图片扩展画面内容。这个视频是介绍如何使用stable diffusion中的outpainting（局部重绘）功能来补充图片以外画面，结合PS的粗略处理，可以得到一个完美画面。让AI成为画手的一个得力工具。, 视频播放量 14221、弹幕. Apple event, protože nějaký teď nedávno byl. Stable diffusion is an open-source technology. . Depending on how stable diffusion works, it might be interesting to use it to generate. 이제 부터 Stable Diffusion은 줄여서 SD로 표기하겠습니다. Get an approximate text prompt, with style, matching an image. PromptMateIO • 7 mo. safetensor and install it in your "stable-diffusion-webuimodelsStable-diffusion" directory. flickr30k. Then, run the model: import Replicate from "replicate"; const replicate = new Replicate( { auth: process. ckpt checkpoint was downloaded), run the following: Technical details regarding Stable Diffusion samplers, confirmed by Katherine: - DDIM and PLMS are originally the Latent Diffusion repo DDIM was implemented by CompVis group and was default (slightly different update rule than the samplers below, eqn 15 in DDIM paper is the update rule vs solving eqn 14's ODE directly) While Stable Diffusion doesn't have a native Image-Variation task, the authors recreated the effects of their Image-Variation script using the Stable Diffusion v1-4 checkpoint. ComfyUI seems to work with the stable-diffusion-xl-base-0. Write a logo prompt and watch as the A. The weights were ported from the original implementation. On SD 2. Model Overview. 手順3：学習を行う. A text-guided inpainting model, finetuned from SD 2. 160 upvotes · 39 comments. Check it out: Stable Diffusion Photoshop Plugin (0. Uses pixray to generate an image from text prompt. 0 和 2. Mage Space and Yodayo are my recommendations if you want apps with more social features. StableDiffusion. Stable Diffusion. 9 on ubuntu 22. Img2txt. Apply the filter: Apply the stable diffusion filter to your image and observe the results. Bootstrapping Language-Image Pre-training. 主にテキスト入力に基づく画像生成（text-to-image）に使用されるが、他にも インペインティング （ 英語版. Doing this on a loop takes advantage of the imprecision in using CLIP latent space walk - fixed seed but two different prompts. 002. 0 was released in November 2022 and has been entirely funded and developed by Stability AI. 今回つくった画像はこんなのになり. Prompt string along with the model and seed number. Flirty_Dane • 7 mo. (You can also experiment with other models. Using the above metrics helps evaluate models that are class-conditioned. jpeg by default on the root of the repo. 6 The Stable Diffusion 2 repository implemented all the servers in gradio and streamlit model-type is the type of image modification demo to launch For example, to launch the streamlit version of the image upscaler on the model created in the original step (assuming the x4-upscaler-ema. Want to see examples of what you can build with Replicate? Check out our showcase. The results from the Stable Diffusion and Kandinsky models vary due to their architecture differences and training process; you can generally expect SDXL to produce higher quality images than Stable Diffusion v1. stable-diffusion-LOGO-fine-tuned model trained by nicky007. Stable Diffusion （ステイブル・ディフュージョン）は、2022年に公開された ディープラーニング （深層学習）の text-to-imageモデル （ 英語版 ） である。. Option 2: Install the extension stable-diffusion-webui-state. Check out the Quick Start Guide if you are new to Stable Diffusion. 9 fine, but when I try to add in the stable-diffusion. 手順1：教師データ等を準備する. However, at the time he installed it only one . September 14, 2022 AI/ML. photo of perfect green apple with stem, water droplets, dramatic lighting. This version is optimized for 8gb of VRAM. stable-diffusion. ago Stable diffusion uses openai clip for img2txt and it works pretty well. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. The StableDiffusionImg2ImgPipeline uses the diffusion-denoising mechanism proposed in SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations by Chenlin. Stable Diffusion XL (SDXL) Inpainting. In your stable-diffusion-webui folder, create a sub-folder called hypernetworks. Others are delightfully strange. Další příspěvky na téma Stable Diffusion. ai and more. . LoRAモデルを使って画像を生成する方法（Stable Diffusion web UIが必要）. The vulnerability has been addressed in Ghostscript 9. Txt2Img：文生图 Img2Txt：图生文 Img2Img：图生图 功能点 部署 Stable Diffusion WebUI 更新 python 版本 切换国内 Linux 安装镜像 安装 Nvidia 驱动 安装stable-diffusion-webui 并启动服务 部署飞书机器人 操作方式 操作命令 设置关键词: 探索企联AI Hypernetworks. So 4 seeds per prompt, 8 total. But the […]Stable Diffusion是2022年發布的深度學習 文本到图像生成模型。 它主要用於根據文本的描述產生詳細圖像，儘管它也可以應用於其他任務，如內補繪製、外補繪製，以及在提示詞指導下產生圖生圖的转变。. Controlnet面部控制，完美复刻人脸 (基于SD2. Stable Diffusion XL. ArtBot or Stable UI are completely free, and let you use more advanced Stable Diffusion features (such as. Lexica is a collection of images with prompts. For more details on how this dataset was scraped, see Midjourney User. Run time and cost. Navigate to txt2img tab, find Amazon SageMaker Inference panel. At the time of release (October 2022), it was a massive improvement over other anime models. Just two. Qualcomm has demoed AI image generator Stable Diffusion running locally on a mobile in under 15 seconds. Mac: run the command . 前提：Stable. . More awesome work from Christian Cantrell in his free plugin. It is defined simply as a dilation followed by an erosion using the same structuring element used in the opening operation. Enter a prompt, and click generate. 0 release includes robust text-to-image models trained using a brand new text encoder (OpenCLIP), developed by LAION with support. The extensive list of features it offers can be intimidating. 1 Model Cards (768x768px) - Model Cards/Weights for Stable Diffusion 2. This extension adds a tab for CLIP Interrogator. Stability AI는 방글라데시계 영국인. This is no longer the case. 81 seconds. . About that huge long negative prompt list. Appendix A: Stable Diffusion Prompt Guide. (Optimized for stable-diffusion (clip ViT-L/14)) Public; 2. 8 pip install torch torchvision -. with current technology would it be possible to ask the AI to generate a text from an image? in order to know what technology could describe the image, a tool for AI to describe the image for us. The inspiration was simply the lack of any Emiru model of any sort here. Stable diffusion image-to-text (SDIT) is an advanced image captioning model based on the GPT architecture and uses a diffusion-based training algorithm to improve stability and consistency during training. Dear friends, come and join me on an incredible journey through Stable Diffusion. We assume that you have a high-level understanding of the Stable Diffusion model. Step 3: Clone web-ui. Text-to-image. The Stable Diffusion 2. Only text prompts are provided. Copy the prompt, paste it to the Stable Diffusion and press Generate to see generated images. r/StableDiffusion. 1 I use this = oversaturated, ugly, 3d, render, cartoon, grain, low-res, kitsch, black and white. If you are absolutely sure that the AI image you want to extract the prompt from was generated using Stable Diffusion, then this method is just for you. To quickly summarize: Stable Diffusion (Latent Diffusion Model) conducts the diffusion process in the latent space, and thus it is much faster than a pure diffusion model. 0 (SDXL 1. Just go to this address and you will see and learn: Fine-tune Your AI Images With These Simple Prompting Techniques - Stable Diffusion Art (stable-diffusion-art. 5. This guide will show you how to finetune DreamBooth. The average face of a teacher generated by Stable Diffusion and DALL-E 2. For 2. Running App Files Files Community 37 Discover amazing ML apps made by the community. A surrealist painting of a cat by Salvador Dali/r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. This model can follow a two-stage model process (though each model can also be used alone); the base model generates an image, and a refiner model takes that image and further enhances its details and quality. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. To use this pipeline for image-to-image, you’ll need to prepare an initial image to pass to the pipeline. A fun little AI art widget named Text-to-Pokémon lets you plug in any name or. More info: Discord: Check out our new Lemmy instance. Get inspired with Kiwi Prompt's stable diffusion prompts for clothes. safetensors format. Then, select the base image and additional references for details and styles. We build on top of the fine-tuning script provided by Hugging Face here. Settings for all eight stayed the same: Steps: 20, Sampler: Euler a, CFG scale: 7, Face restoration: CodeFormer, Size: 512x768, Model hash: 7460a6fa. stable-diffusion txt2img参数整理 Sampling steps ：采样步骤”：“迭代改进生成图像的次数；较高的值需要更长的时间；非常低的值可能会产生糟糕的结果”， 指的是Stable Diffusion生成图像所需的迭代步数。Stable Diffusion is a cutting-edge text-to-image diffusion model that can generate photo-realistic images based on any given text input. First, your text prompt gets projected into a latent vector space by the. img2txt OR "prompting" is the reverse operation, convergent, from significantly many more bits to significantly less or small count of bits, like a capture card does, but. 1 1 comment Evnl2020 • 1 yr. It’s a fun and creative way to give a unique twist to my images. Para hacerlo, tienes que registrarte en la web beta. Credit Calculator. A decoder, which turns the final 64x64 latent patch into a higher-resolution 512x512 image. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. This distribution is changing rapidly. Public. Stable diffusion image-to-text (SDIT) is an advanced image captioning model based on the GPT architecture and uses a diffusion-based training algorithm to improve stability and. Given a (potentially crude) image and the right text prompt, latent diffusion. ckpt for using v1. Inside your subject folder, create yet another subfolder and call it output. To put another way, quoting your source at gigazine, "the larger the CFG scale, the more likely it is that a new image can be generated according to the image input by the prompt. Go to extensions tab; Click "Install from URL" sub tabtry going to an image editor like photoshop or gimp, find a picture of crumpled up paper, something that has some textures in it and use it as a background, add your logo on the top layer and apply some small amount of noise to the whole thing, make sure to have a good amount of contrast between the background and foreground (if your background. On the first run, the WebUI will download and install some additional modules. This version is optimized for 8gb of VRAM. 667 messages. During our research, jp2a , which works similarly to img2txt, also appeared on the scene. It’s a simple and straightforward process that doesn’t require any technical expertise. comments sorted by Best Top New Controversial Q&A Add a Comment. ) Come up with a prompt that describe your final picture as accurately as possible. 2. card. Para ello vam. Hieronymus Bosch. Stable Diffusion without UI or tricks (only take off filter xD). I was using one but it does not work anymore since yesterday. Updated 1 day, 17 hours ago 53 runs fofr / sdxl-pixar-cars SDXL fine-tuned on Pixar Cars. Unprompted is a highly modular extension for AUTOMATIC1111's Stable Diffusion Web UI that allows you to include various shortcodes in your prompts. A method to fine tune weights for CLIP and Unet, the language model and the actual image de-noiser used by Stable Diffusion, generously donated to the world by our friends at Novel AI in autumn 2022. You can also upload and replicate non-AI generated images. Additional training is achieved by training a base model with an additional dataset you are. You are welcome to try our free online Stable Diffusion based image generator at It supports img2img generation, including sketching of the initial image :) Cool site. Discover stable diffusion Img2Img techniques & their applications. It uses the Stable Diffusion x4 upscaler. GitHub. 画像から画像を作成する. StableDiffusion - Txt2Img - HouseofCat Stable Diffusion 2. Fine-tuned Model Checkpoints (Dreambooth Models) Download the custom model in Checkpoint format (. ChatGPT page. Start the WebUI. All you need to do is to download the embedding file stable-diffusion-webui > embeddings and use the extra. 1M runs. [1] Generated images are. 13:23. 0 前回 1. 7>"), and on the script's X value write something like "-01, -02, -03", etc. a. lupaspirit. txt2img OR "imaging" is mathematically divergent operation, from less bits to more bits, even ARM or RISC-V can do that. Stable Diffusion. 1. dreamstudio. Get prompts from stable diffusion generated images. Explore and run machine. 2. This example was created by a different version, rmokady/clip_prefix_caption:d703881e. Our AI-generated prompts can help you come up with. 0) Watch on. Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. 0. ago. If the image with the text was clear enough, you will receive recognized and readable text. Easy Prompt SelectorのYAMLファイルは「stable-diffusion-webuiextensionssdweb-easy-prompt-selector ags」の中にあります。 「. true. k. Stable DiffusionはNovelAIやMidjourneyとはどう違うの？ Stable Diffusionを簡単に使えるツールは結局どれを使えばいいの？ 画像生成用のグラフィックボードを買うならどれがオススメ？ モデルのckptとsafetensorsって何が違うの？ モデルのfp16・fp32・prunedって何？本教程需要一些AI绘画基础，并不是面对0基础人员，如果你没有学习过stable diffusion的基本操作或者对Controlnet插件毫无了解，可以先看看秋葉aaaki等up的教程，做到会存放大模型，会安装插件并且有基本的视频剪辑能力。-----一、准备工作This issue is a workaround for a security vulnerability. like 4. Output. Press the Window key (It should be on the left of the space bar on your keyboard), and a search window should appear. 缺點：. The pre-training dataset of Stable Diffusion may have limited overlap with the pre-training dataset of InceptionNet, so it is not a good candidate here for feature extraction. Example outputs . 使用anaconda进行webui的创建. ·. AIイラストに衣装を着せたときの衣装の状態に関する呪文（プロンプト）についてまとめました。 七海が実際にStable Diffusionで生成したキャラクターを使って検証した衣装の状態に関する呪文をご紹介します。 ※このページから初めて、SThis tutorial shows how to fine-tune a Stable Diffusion model on a custom dataset of {image, caption} pairs. Generate the image. Does anyone know of any extensions for A1111, that allow you to insert a picture, and it can give you a prompt? I tried a feature like it on my. It is our fastest API, matching the speed of its predecessor, while providing higher quality image generations at 512x512 resolution. A k tomu “man struck down” kde už vlastně ani nevím proč jsem to potřeboval. • 7 mo. Linux: run the command webui-user. 上个月做了安卓和苹果手机用远端sd进行跑图的几个demo，整体流程很简单. ckpt (5. ← Runway previews text to video Lexica: Search for AI-made art, with prompts →. This may take a few minutes. Using a model is an easy way to achieve a certain style. To use a VAE in AUTOMATIC1111 GUI, go to the Settings tab and click the Stabe Diffusion section on the left. Enter the required parameters for inference. In this quick episode we do a simple workflow where we upload an image into our SDXL graph inside of ComfyUI and add additional noise to produce an altered i. Text-to-image models like Stable Diffusion generate an image from a text prompt. 5、2. On Ubuntu 19. If you have 8gb RAM, consider making an 8gb page file/swap file, or use the --lowram option (if you have more gpu vram than ram). 指定した画像に近づくように画像生成する機能です。通常のプロンプトによる生成指定に加えて、追加でVGG16の特徴量を取得し、生成中の画像が指定したガイド画像に近づくよう、生成される画像をコントロールします。 2. Roboti na kole. 多種多様な表現が簡単な指示で行えるようになり、人間の負担が著しく減ります。. fixは高解像度の画像が生成できるオプションです。. A decoder, which turns the final 64x64 latent patch into a higher-resolution 512x512 image. Stable Horde for Web UI. • 7 mo. they converted to a. Negative embeddings bad artist and bad prompt. The StableDiffusionImg2ImgPipeline uses the diffusion-denoising mechanism proposed in SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations by. Press “+ New Chat” button on the left panel to start a new conversation. C:stable-diffusion-uimodelsstable-diffusion)Option 1: Every time you generate an image, this text block is generated below your image. Stable Diffusion img2img support comes to Photoshop. 0) のインストール，画像生成（img2txt），画像変換（img2img），APIを利用して複数画像を一括生成（AUTOMATIC1111，Python，PyTorch を使用）（Windows 上）Step#1: Setup your environment. Use the resulting prompts with text-to-image models like Stable Diffusion to create cool art! Public. 尚未安裝 Stable Diffusion WebUI 的夥伴可以參考上一篇 如何在 M1 Macbook 上跑 Stable Diffusion？Stable Diffusion Checkpoint: Select the model you want to use. More posts you may like r/selfhosted Join • 13. 0 的过程，包括下载必要的模型以及如何将它们安装到. 2. What is Img2Img in Stable Diffusion Setting up The Software for Stable Diffusion Img2img How to Use img2img in Stable Diffusion Step 1: Set the. 0) Watch on. 5 it/s. I do think that your approach will struggle by the fact it's a similar training method on the already limited faceset you have - so if it's not good enough to work already in DFL for producing those missing angles I'm not sure stable-diffusion will let you. Next and SDXL tips. Stable diffusion is a critical aspect of obtaining high-quality image transformations using Img2Img. 5. Payload is a config-based, code-first CMS and application framework. . img2txt arch. Two main ways to train models: (1) Dreambooth and (2) embedding. We recommend to explore different hyperparameters to get the best results on your dataset. The model bridges the gap between vision and natural. I have been using Stable Diffusion for about 2 weeks now. 4); stable_diffusion (v1. Write a logo prompt and watch as the A. 1M runs. 4 ・diffusers 0. It scaffolds the data that Payload stores as well as maintains custom React components, hook logic, custom validations, and much more. By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. The comparison of SDXL 0. Go to extensions tab; Click "Install from URL" sub tab try going to an image editor like photoshop or gimp, find a picture of crumpled up paper, something that has some textures in it and use it as a background, add your logo on the top layer and apply some small amount of noise to the whole thing, make sure to have a good amount of contrast between the background and foreground (if your background. plugin already! NOTE: Once installed, you will be able to generate images without a subscrip. Running the Diffusion Process. File "scriptsimg2txt. 5, Stable Diffusion XL (SDXL), and Kandinsky 2. 04 through 22. r/StableDiffusion •. . com. Negative prompting influences the generation process by acting as a high-dimension anchor,. The program needs 16gb of regular RAM to run smoothly. Drag and drop an image image here (webp not supported). Fix it to look like the original. It's stayed fairly consistent with Img2Img batch processing. py", line 144, in interrogate load_blip_model(). It. Image: The Verge via Lexica. r/sdnsfw Lounge. Predictions typically complete within 27 seconds. Take the “Behind the scenes of the moon landing” image. Stable Diffusion lets you create images using just text prompts but if you want them to look stunning, you must take advantage of negative prompts. Uncrop. I wanted to report some observations and wondered if the community might be able to shed some light on the findings. Text prompt with description of the things you want in the image to be generated. • 1 yr. With its 860M UNet and 123M text encoder. Get an approximate text prompt, with style, matching an image. Embeddings (aka textual inversion) are specially trained keywords to enhance images generated using Stable Diffusion. But the width, height and other defaults need changing. 手順3：学習を行う. Introduction; Architecture; RequirementThe Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images. Use. img2txt github. Put the Lora of the first epoch in your prompt (like "<lora:projectname-01:0. The default value is set to 2. ai, y. Functioning as image viewers for the terminal, chafa and catimg have only been an integral part of a stable release of the Linux distribution since Debian GNU/Linux 10. The text to image sampling script within Stable Diffusion, known as "txt2img", consumes a text prompt in addition to assorted option parameters covering sampling types, output image dimensions, and seed values. You are welcome to try our free online Stable Diffusion based image generator at It supports img2img generation, including sketching of the initial image :) Cool site. Discover amazing ML apps made by the communityPosition the 'Generation Frame' in the right place.